# Creative Design

This is a DreamBooth LoRA model based on HiDream-ai/HiDream-I1-Full, specializing in generating yarn art-style dog images.

An unconditional image generation model based on the diffusers library, focusing on animal image generation.

A text-to-image generation model based on black-forest-labs/FLUX.1-dev foundation model, fine-tuned using Low-Rank Adaptation (LoRA) technology

doa is a text-to-image generation project created based on the SDXL model, developed by reijlita, which can achieve high-quality text-to-image conversion.

FLUX.1-dev is a text-to-image generation model based on the diffusers library, focusing on high-quality image synthesis.

A text-to-image generation model based on StableDiffusionXLPipeline, focused on artistic creation.

SDXL-Turbo is an efficient text-to-image generation model capable of producing high-quality images with very few sampling steps.

PixelGen is a text-to-image generation model developed by HelpingAI, capable of producing high-quality images from textual descriptions.

This is a Low-Rank Adaptation (LoRA) model based on Stable Diffusion XL, specifically designed to generate highly aesthetic images.

This is a fine-tuned model based on Stable Diffusion XL, specializing in generating eye-catching style images.

This is a LoRA model based on Stable Diffusion XL (SDXL), specifically designed for generating paper cutout-style images.

This model can generate high-quality images based on input text descriptions, suitable for various scenarios such as creative design and content creation.

This is an unconditional image generation model based on the Apache-2.0 license, primarily used for image generation tasks in the field of art creation.

PortraitPlus is a text-to-image generation model based on Stable Diffusion, specializing in generating high-quality portraits.

This is a text-to-image generation model fine-tuned based on Stable Diffusion v1-5, specifically designed for generating emoji-style images.